Last Updated: November 21, 2025
Stable Diffusion
AI image generation
Core Concepts
| Item | Description |
|---|---|
Text Prompt
|
Description of desired image |
Negative Prompt
|
What to avoid |
Steps
|
Denoising iterations (20-50) |
CFG Scale
|
Prompt adherence (7-15) |
Sampler
|
Denoising algorithm |
Seed
|
Reproducibility number |
Using Diffusers Library
from diffusers import StableDiffusionPipeline
import torch
# Load model
pipe = StableDiffusionPipeline.from_pretrained(
"stabilityai/stable-diffusion-2-1",
torch_dtype=torch.float16
).to("cuda")
# Generate image
image = pipe(
prompt="a cat wearing sunglasses, digital art",
negative_prompt="blurry, low quality",
num_inference_steps=50,
guidance_scale=7.5,
height=512,
width=512
).images[0]
image.save("output.png")
Prompt Tips
- Be specific and descriptive
- Add style keywords (oil painting, 3D render, etc.)
- Use negative prompts to avoid unwanted elements
- Include quality modifiers (highly detailed, 4k, etc.)
- Reference artists or art styles
- Adjust CFG scale if image doesn't match prompt
Best Practices
- Start with 512x512 for speed
- Use 20-50 steps (more steps = diminishing returns)
- Save seed for reproducibility
- Use img2img for variations
💡 Pro Tips
Quick Reference
Lower CFG scale for more creative, higher for more literal